Individual skill

stable-diffusion-image-generation

Generates images with Stable Diffusion models (SD 1.5, SDXL, SD 3.0, Flux) through the HuggingFace Diffusers library, covering text-to-image, image-to-image, inpainting, outpainting, ControlNet conditioning, LoRA adapters, scheduler swapping, and GPU memory optimization. Use when generating images from text prompts, transforming or restyling an existing image, filling masked regions of an image, adding spatial control from edges, poses, or depth maps, loading LoRA style adapters, or fixing out-of-memory and black-image errors in a Diffusers pipeline. Not for API-only generation without a GPU, where DALL-E 3 fits better.

plugins/KalarisLabs/research-agent-skills/skills/stable-diffusion-image-generation/SKILL.md

Source commit d780783eba4ecc7c8b9429f95c9362df95f07c61

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Yes

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No

Observed behavior

Network or broad shell

Last checked

Nov 4, 2026

These are observations at the recorded commit, not a guarantee of safety. Review the source before running it.

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